My own implementation of Stable Diffusion for me to generate reference art

Image generation by text prompts, using Stable Diffusion-v1-5, in PyTorch & Gradio

Easy to use diffusion in mobile interface.

Using SageMaker and LoRA to fine-tune the Stable Diffusion model and generate fashion images

한국어 특화 이미지 생성 (롯데 이노베이트 인턴십)

Generate image by prompts using Stable Diffusion XL-turbo, in Pytorch & Gradio

Generate photo-realistic & high-resolution images by user-defined prompts using Flux-schnell, in PyTorch & Gradio

Here I have used Stable Diffusion with the diffusers from hugging face, took one image as input and then added elements into it by using diffusion algorithm, and iterated this process three/four times with different type of elements, adding by text inputs.

This is an adaption of the notebook , which is provided as part of a class by Hugging Face

Generating images with diffusion models on a mobile device, with an intranet GPU box as backend

a tool to create synthetic image data

This project is a Streamlit-based web app that enables users to generate AI-generated images using text prompts. The app integrates two powerful image generation models: OpenAI's DALL-E and Huggingface's Diffusion models.

Experimental Stable Diffusion XL Webui

Configuration files for building E621-Rising v3 SDXL model and dataset

A web app that allows you to select a subject and then change its background, OR keep the background and change the subject.

diffusion model for unconditional image generation of Bored Apes

This is a repo providing same stable diffusion experiments, regarding textual inversion task and captioning task

Implementation of Paint-with-words with Stable Diffusion using diffusers pipeline: method from eDiff-I that let you generate image from text-labeled segmentation map.

Easily create your own AI avatar images!

Toolchain for creating custom datasets and training Stable Diffusion (1.x, 2.x, XL) models and LoRAs